Flux.1 Image Generation

Unleash your creativity with Flux.1, a powerful diffusion model for text-to-image generation and advanced image editing.

Platform: Replicate

Image GenerationImage InpaintingLoRA AdaptationDiffusion Model

32 runs

H100

License Check Required

🚀Function Overview

A diffusion-based model for text-to-image generation, image-to-image transformation, and masked inpainting with adjustable parameters for resolution, style intensity, and LoRA integration.

Key Features

Generates images from text prompts with optional trigger words
Supports image-to-image transformation and inpainting via masks
Customizable aspect ratios, resolutions, and output formats
Adjustable denoising steps, guidance scale, and prompt strength
LoRA model integration for style/concept adaptation with scaling controls
FP8 quantization for fast generation mode

Use Cases

•Creative artwork generation from text descriptions
•Photo editing and object removal through inpainting
•Style transfer using custom LoRA models
•Rapid prototyping of visual concepts

⚙️Input Parameters

prompt

string

Prompt for generated image. If you include the `trigger_word` used in the training process you are more likely to activate the trained object, style, or concept in the resulting image.

image

string

Input image for image to image or inpainting mode. If provided, aspect_ratio, width, and height inputs are ignored.

mask

string

Image mask for image inpainting mode. If provided, aspect_ratio, width, and height inputs are ignored.

aspect_ratio

string

Aspect ratio for the generated image. If custom is selected, uses height and width below & will run in bf16 mode

height

integer

Height of generated image. Only works if `aspect_ratio` is set to custom. Will be rounded to nearest multiple of 16. Incompatible with fast generation

width

integer

Width of generated image. Only works if `aspect_ratio` is set to custom. Will be rounded to nearest multiple of 16. Incompatible with fast generation

prompt_strength

number

Prompt strength when using img2img. 1.0 corresponds to full destruction of information in image

model

string

Which model to run inference with. The dev model performs best with around 28 inference steps but the schnell model only needs 4 steps.

num_outputs

integer

Number of outputs to generate

num_inference_steps

integer

Number of denoising steps. More steps can give more detailed images, but take longer.

guidance_scale

number

Guidance scale for the diffusion process. Lower values can give more realistic images. Good values to try are 2, 2.5, 3 and 3.5

seed

integer

Random seed. Set for reproducible generation

output_format

string

Format of the output images

output_quality

integer

Quality when saving the output images, from 0 to 100. 100 is best quality, 0 is lowest quality. Not relevant for .png outputs

disable_safety_checker

boolean

Disable safety checker for generated images.

go_fast

boolean

Run faster predictions with model optimized for speed (currently fp8 quantized); disable to run in original bf16

megapixels

string

Approximate number of megapixels for generated image

lora_scale

number

Determines how strongly the main LoRA should be applied. Sane results between 0 and 1 for base inference. For go_fast we apply a 1.5x multiplier to this value; we've generally seen good performance when scaling the base value by that amount. You may still need to experiment to find the best value for your particular lora.

extra_lora

string

Load LoRA weights. Supports Replicate models in the format <owner>/<username> or <owner>/<username>/<version>, HuggingFace URLs in the format huggingface.co/<owner>/<model-name>, CivitAI URLs in the format civitai.com/models/<id>[/<model-name>], or arbitrary .safetensors URLs from the Internet. For example, 'fofr/flux-pixar-cars'

extra_lora_scale

number

Determines how strongly the extra LoRA should be applied. Sane results between 0 and 1 for base inference. For go_fast we apply a 1.5x multiplier to this value; we've generally seen good performance when scaling the base value by that amount. You may still need to experiment to find the best value for your particular lora.

💡Usage Examples

Example 1

Input Parameters

{
  "model": "dev",
  "prompt": "TOKLEE wearing a black tracksuit and black cap in new york mixing studio laying vocals in the booth headphones and reading lyrics from  a piece of paper, TOKLEE is standing and behind the glass, we see the engineer from behind the mixing desk then TOKLEE face on behind the glass of recording booth",
  "go_fast": false,
  "lora_scale": 1,
  "megapixels": "1",
  "num_outputs": 1,
  "aspect_ratio": "1:1",
  "output_format": "webp",
  "guidance_scale": 1.97,
  "output_quality": 80,
  "prompt_strength": 0.8,
  "extra_lora_scale": 1,
  "num_inference_steps": 28
}

Output Results

https://replicate.delivery/xezq/el3ZGBmeY4sdZkisKLO4il04OKeOFAXqlBHEh5pUGuVeeZalC/out-0.webp

Quick Actions

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Technical Specifications

Hardware Type: H100
Run Count: 32
Commercial Use: Unknown/Restricted
Platform: Replicate

Related Keywords

Text-to-Image GenerationImage InpaintingLoRA IntegrationDiffusion ModelImage-to-Image TransformationResolution ControlCreative Artwork GenerationPhoto Editing

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