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Menilub AI Image Generator

Discover Menilub, a powerful AI model for sophisticated image generation and editing. Ready to experience the power of AI? Start your journey here!

Platform: Replicate
Text-to-Image DiffusionImage InpaintingLoRA IntegrationImage-to-Image Translation
6 runs
H100
License Check Required

🚀Function Overview

A diffusion-based model that generates customized images from text prompts, performs inpainting with masks, and allows style control via LoRA weights, with adjustable parameters for resolution, quality, and generation speed.

Key Features

  • Text-to-image generation with detailed prompt control
  • Image inpainting and image-to-image transformation
  • LoRA weight integration for style/concept tuning
  • Multiple resolution/aspect ratio options
  • Configurable denoising steps and guidance scale
  • Fast generation mode (fp8 quantized)
  • Output quality and format customization

Use Cases

  • Creating custom artwork from text descriptions
  • Photo editing and object removal via inpainting
  • Style transfer using LoRA weights
  • Rapid prototyping of visual concepts
  • Generating variations of existing images

⚙️Input Parameters

prompt

string

Prompt for generated image. If you include the `trigger_word` used in the training process you are more likely to activate the trained object, style, or concept in the resulting image.

image

string

Input image for image to image or inpainting mode. If provided, aspect_ratio, width, and height inputs are ignored.

mask

string

Image mask for image inpainting mode. If provided, aspect_ratio, width, and height inputs are ignored.

aspect_ratio

string

Aspect ratio for the generated image. If custom is selected, uses height and width below & will run in bf16 mode

height

integer

Height of generated image. Only works if `aspect_ratio` is set to custom. Will be rounded to nearest multiple of 16. Incompatible with fast generation

width

integer

Width of generated image. Only works if `aspect_ratio` is set to custom. Will be rounded to nearest multiple of 16. Incompatible with fast generation

prompt_strength

number

Prompt strength when using img2img. 1.0 corresponds to full destruction of information in image

model

string

Which model to run inference with. The dev model performs best with around 28 inference steps but the schnell model only needs 4 steps.

num_outputs

integer

Number of outputs to generate

num_inference_steps

integer

Number of denoising steps. More steps can give more detailed images, but take longer.

guidance_scale

number

Guidance scale for the diffusion process. Lower values can give more realistic images. Good values to try are 2, 2.5, 3 and 3.5

seed

integer

Random seed. Set for reproducible generation

output_format

string

Format of the output images

output_quality

integer

Quality when saving the output images, from 0 to 100. 100 is best quality, 0 is lowest quality. Not relevant for .png outputs

disable_safety_checker

boolean

Disable safety checker for generated images.

go_fast

boolean

Run faster predictions with model optimized for speed (currently fp8 quantized); disable to run in original bf16

megapixels

string

Approximate number of megapixels for generated image

lora_scale

number

Determines how strongly the main LoRA should be applied. Sane results between 0 and 1 for base inference. For go_fast we apply a 1.5x multiplier to this value; we've generally seen good performance when scaling the base value by that amount. You may still need to experiment to find the best value for your particular lora.

extra_lora

string

Load LoRA weights. Supports Replicate models in the format <owner>/<username> or <owner>/<username>/<version>, HuggingFace URLs in the format huggingface.co/<owner>/<model-name>, CivitAI URLs in the format civitai.com/models/<id>[/<model-name>], or arbitrary .safetensors URLs from the Internet. For example, 'fofr/flux-pixar-cars'

extra_lora_scale

number

Determines how strongly the extra LoRA should be applied. Sane results between 0 and 1 for base inference. For go_fast we apply a 1.5x multiplier to this value; we've generally seen good performance when scaling the base value by that amount. You may still need to experiment to find the best value for your particular lora.

💡Usage Examples

Example 1

Input Parameters

{
  "model": "dev",
  "prompt": "Dynamic cartoon illustration: 56-year-old male Menilub, depicted as a friendly software development team leader, enthusiastically explaining a complex idea on a giant whiteboard filled with quirky diagrams, arrows, and funny code snippets. He has slightly oversized, expressive eyes (no glasses), kind smile lines, very short silvering hair, and wears a smart-casual collared shirt. His team (simplified cartoon characters) look on with interest. Bright, clean office background. Focus on Menilub's clear, animated, and unobstructed face, capturing his distinct friendly yet leader-like look. Style: Modern, clean-line cartoon.",
  "go_fast": false,
  "lora_scale": 1,
  "megapixels": "1",
  "num_outputs": 1,
  "aspect_ratio": "1:1",
  "output_format": "webp",
  "guidance_scale": 3,
  "output_quality": 80,
  "prompt_strength": 0.8,
  "extra_lora_scale": 1,
  "num_inference_steps": 28
}

Output Results

https://replicate.delivery/xezq/bqG4Sfv7p8W2YaK6wf05FXTBORRJmdwGy5lC8Qz0v6ZfIKNpA/out-0.webp

Quick Actions

Technical Specifications

Hardware Type
H100
Run Count
6
Commercial Use
Unknown/Restricted
Platform
Replicate

Related Keywords

Text-to-image generationImage inpaintingLoRA integrationPhoto editingStyle transferRapid visual prototypingCustom artworkImage variations