QuiGon Live 2 Image Generation

Discover QuiGon Live 2 Image Generation, a powerful model for advanced image creation and editing. Let's explore what this AI model can do for you!

Platform: Replicate

Image GenerationImage InpaintingLoRA IntegrationCustomizable Diffusion

71 runs

H100

License Check Required

🚀Function Overview

Generates and edits images using diffusion-based techniques with configurable parameters, supporting text prompts, image-to-image transformations, mask-based inpainting, LoRA model integration, and output quality controls.

Key Features

Prompt-based image generation with trigger word optimization
Image-to-image transformation and inpainting via masks
Aspect ratio and resolution controls (with bf16/fp8 modes)
Choice between quality-optimized ('dev') or speed-optimized ('schnell') models
Adjustable denoising steps and guidance scale
External LoRA model integration from Replicate/HuggingFace/CivitAI
Output format/quality customization and safety checker toggle

Use Cases

•Creating character art from descriptive prompts
•Editing existing images through inpainting
•Style transfer using custom LoRA models
•Batch image generation with precise resolution control

⚙️Input Parameters

prompt

string

Prompt for generated image. If you include the `trigger_word` used in the training process you are more likely to activate the trained object, style, or concept in the resulting image.

image

string

Input image for image to image or inpainting mode. If provided, aspect_ratio, width, and height inputs are ignored.

mask

string

Image mask for image inpainting mode. If provided, aspect_ratio, width, and height inputs are ignored.

aspect_ratio

string

Aspect ratio for the generated image. If custom is selected, uses height and width below & will run in bf16 mode

height

integer

Height of generated image. Only works if `aspect_ratio` is set to custom. Will be rounded to nearest multiple of 16. Incompatible with fast generation

width

integer

Width of generated image. Only works if `aspect_ratio` is set to custom. Will be rounded to nearest multiple of 16. Incompatible with fast generation

prompt_strength

number

Prompt strength when using img2img. 1.0 corresponds to full destruction of information in image

model

string

Which model to run inference with. The dev model performs best with around 28 inference steps but the schnell model only needs 4 steps.

num_outputs

integer

Number of outputs to generate

num_inference_steps

integer

Number of denoising steps. More steps can give more detailed images, but take longer.

guidance_scale

number

Guidance scale for the diffusion process. Lower values can give more realistic images. Good values to try are 2, 2.5, 3 and 3.5

seed

integer

Random seed. Set for reproducible generation

output_format

string

Format of the output images

output_quality

integer

Quality when saving the output images, from 0 to 100. 100 is best quality, 0 is lowest quality. Not relevant for .png outputs

disable_safety_checker

boolean

Disable safety checker for generated images.

go_fast

boolean

Run faster predictions with model optimized for speed (currently fp8 quantized); disable to run in original bf16

megapixels

string

Approximate number of megapixels for generated image

lora_scale

number

Determines how strongly the main LoRA should be applied. Sane results between 0 and 1 for base inference. For go_fast we apply a 1.5x multiplier to this value; we've generally seen good performance when scaling the base value by that amount. You may still need to experiment to find the best value for your particular lora.

extra_lora

string

Load LoRA weights. Supports Replicate models in the format <owner>/<username> or <owner>/<username>/<version>, HuggingFace URLs in the format huggingface.co/<owner>/<model-name>, CivitAI URLs in the format civitai.com/models/<id>[/<model-name>], or arbitrary .safetensors URLs from the Internet. For example, 'fofr/flux-pixar-cars'

extra_lora_scale

number

Determines how strongly the extra LoRA should be applied. Sane results between 0 and 1 for base inference. For go_fast we apply a 1.5x multiplier to this value; we've generally seen good performance when scaling the base value by that amount. You may still need to experiment to find the best value for your particular lora.

💡Usage Examples

Example 1

Input Parameters

{
  "model": "dev",
  "prompt": "Qui-Gon Jinn looking slightly to the left, his expression calm yet intense, filled with quiet wisdom and conviction. His long brown hair frames his weathered face, and his piercing eyes reflect both serenity and defiance. Soft, natural light illuminates his features, highlighting the lines of experience and compassion. His Jedi robes drape over his shoulders and chest, textured and worn, symbolizing years of service and spiritual depth. Behind him, the blurred interior of the Jedi Temple adds a sacred, timeless atmosphere. A sense of quiet strength radiates from him—this is a Jedi who listens to the Force above all else.",
  "go_fast": false,
  "lora_scale": 1,
  "megapixels": "1",
  "num_outputs": 1,
  "aspect_ratio": "21:9",
  "output_format": "jpg",
  "guidance_scale": 3,
  "output_quality": 80,
  "prompt_strength": 0.8,
  "extra_lora_scale": 1,
  "num_inference_steps": 28
}

Output Results

https://replicate.delivery/xezq/41MTV1UJp9piL5YtnzHFsPKhK6f8zOQ6kIlJ5ogfdVVqgicUA/out-0.jpg

Quick Actions

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Technical Specifications

Hardware Type: H100
Run Count: 71
Commercial Use: Unknown/Restricted
Platform: Replicate

Related Keywords

Image EditingLoRA IntegrationText-to-ImageImage-to-ImageInpaintingCustomizable Diffusion

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