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All AI Models

Browse model descriptions, inputs, examples and provider information.

Page 13 of 39 · Models 145–156 of 467

OpenAI GPT-4o-mini

openai/gpt-4o-mini

OpenAI deprecated GPT-4o mini on October 14, 2026, with listed API snapshots due to shut down January 14, 2027. Verify Replicate access separately.

Text GenerationMultimodal AIConversational Assistant+1
1.5M
Priced by multiple properties
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Familia Alves Image Generator

carlaluiza10/familiaalves

Familia Alves Image Generator generates and modifies images based on text prompts and input images, featuring inpainting capabilities, adjustable styles via LoRA, and control over image properties like resolution and aspect ratio.

Image GenerationImage InpaintingLoRA Style Transfer+1
ACE-STEP

lucataco/ace-step

Commercial

ACE-STEP is a diffusion-based foundation model for generating music from text prompts and structured lyrics, optimizing for speed, coherence, and controllability.

Text-to-Music GenerationAI ComposerLyric-Guided Audio
3.6k
L40S
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Sector SmokeBox

baoliwaters/sector_smokebox

Generate and edit images with Sector SmokeBox, supporting text-to-image, inpainting, and LoRA style adjustments for versatile visual content creation.

Image GenerationImage InpaintingText-to-Image+1
Buzz by The AI Surfer Team

theaisurferteam/buzz

Buzz by The AI Surfer Team is a diffusion-based model for generating and editing images through text prompts, supporting inpainting, custom aspect ratios, LoRA scaling, and fast-quantized inference modes.

Text-to-ImageImage InpaintingDiffusion Model+1
Lightricks LTX Video Generation

lightricks/ltx-video-0.9.7

Lightricks LTX Video Generation is a model for generating smooth, videos in real time using text prompts or input images, with granular control over video properties.

Video SynthesisText-to-VideoImage-to-Video+1
1.2k
L40S
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FramePack Video Diffusion

zsxkib/framepack

Commercial

FramePack Video Diffusion is a model that generates videos progressively from an initial image using text prompts, implementing non-F1 FramePack video diffusion with efficient memory management.

Video DiffusionProgressive GenerationImage-to-Video
1.8k
A100 (80GB)
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Carlaluiza Image Generation

carlaluiza10/carlaluiza

Carlaluiza Image Generation generates, edits, or inpaints images based on text prompts and input images, with customization options for resolution, style, and output quality.

Image GenerationImage InpaintingLoRA Adaptation
Psivitoriavidal Image Generation and Editing Model

airlonfilho/psivitoriavidal

Psivitoriavidal Image Generation and Editing Model generates and edits images from text prompts, supporting inpainting, image-to-image transformations, and custom dimensions with LoRA integration for style control.

Image GenerationInpaintingLoRA Integration+1
Forfluxcoed Image Generation

codelace/forfluxcoed

Forfluxcoed Image Generation is a versatile AI for advanced text-to-image and inpainting. Customize visuals with LoRA integration and precise controls.

Diffusion ModelImage InpaintingLoRA Integration+1
DreamO by Bytedance

zsxkib/dream-o

Commercial

DreamO by Bytedance is a unified framework for image customization tasks including identity preservation, virtual try-on, and style transfer using reference images and prompts.

Image CustomizationFace Identity PreservationVirtual Try-On+1
715
A100 (80GB)
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Forflux Image Generation

codelace/forfluxmaylads

Forflux Image Generation generates and manipulates images through text prompts, supports inpainting, aspect ratio control, and LoRA-based style customization.

Image GenerationImage InpaintingLoRA Fine-Tuning+1